Here's something interesting I did few days ago.
- Generated images using mixture of different styles of prompts with SDXL Base Model ( using Diffusers )
- Trained a LoRA with them
- Generated again with this LoRA + Prompts used to generate the training set.
Ended up with results with enhanced effects - glitchier, weirder, high def.
Results => https://imgur.com/gallery/vUobKPK
I’m gonna train another LoRA with these generations and repeat the process obviously!
This is a pretty neat way to bypass the 77 token limit in Diffusers and develop tons of more styles now that I think about it.
You can play around with the LoRA at https://replicate.com/galleri5/nammeh ( GitHub account needed )
Will publish it to CivitAI soon.
Here's something interesting I did few days ago.
- Generated images using mixture of different styles of prompts with SDXL Base Model ( using Diffusers )
- Trained a LoRA with them
- Generated again with this LoRA + Prompts used to generate the training set.
Ended up with results with enhanced effects - glitchier, weirder, high def.
Results => https://imgur.com/gallery/vUobKPK
I’m gonna train another LoRA with these generations and repeat the process obviously!
This is a pretty neat way to bypass the 77 token limit in Diffusers and develop tons of more styles now that I think about it.
You can play around with the LoRA at https://replicate.com/galleri5/nammeh ( GitHub account needed )
Will publish it to CivitAI soon.